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Stable Diffusion Image Generation

Skill Verified Active

State-of-the-art text-to-image generation with Stable Diffusion models via HuggingFace Diffusers. Use when generating images from text prompts, performing image-to-image translation, inpainting, or building custom diffusion pipelines.

Purpose

Generate high-quality images from text prompts, perform image transformations, and build custom diffusion workflows using state-of-the-art models.

Features

  • Text-to-Image generation
  • Image-to-Image translation
  • Inpainting and Outpainting
  • ControlNet and LoRA integration
  • Support for various Stable Diffusion models (SDXL, SD 3.0, Flux)

Use Cases

  • Generating images from detailed text descriptions
  • Transforming existing images with style or content guidance
  • Filling in masked regions of images contextually
  • Building custom image generation pipelines with advanced controls

Non-Goals

  • API-based generation without local GPU
  • Artistic or highly stylized outputs (Midjourney)
  • Cloud-specific integrations (Imagen, Leonardo.ai)

Trust

  • info:Issues AttentionWith 17 open and 4 closed issues in the last 90 days, the maintainer engagement appears moderate but could be improved for faster issue resolution.

Installation

npx skills add davila7/claude-code-templates

Runs the Vercel skills CLI (skills.sh) via npx — needs Node.js locally and at least one installed skills-compatible agent (Claude Code, Cursor, Codex, …). Assumes the repo follows the agentskills.io format.

Quality Score

Verified
95 /100
Analyzed about 17 hours ago

Trust Signals

Last commitabout 19 hours ago
Stars27.2k
LicenseMIT
Status
View Source

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