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Stable Diffusion Image Generation

技能 已验证 活跃

State-of-the-art text-to-image generation with Stable Diffusion models via HuggingFace Diffusers. Use when generating images from text prompts, performing image-to-image translation, inpainting, or building custom diffusion pipelines.

目的

Generate high-quality images from text prompts and other inputs using state-of-the-art diffusion models.

功能

  • Text-to-image generation
  • Image-to-image translation
  • Inpainting and outpainting
  • ControlNet and T2I-Adapter conditioning
  • LoRA adapter support for fine-tuning styles
  • SDXL and SD 3.0 model support
  • Memory optimization techniques
  • Image upscaling and refinement

使用场景

  • Generating artwork from descriptive prompts
  • Transforming existing photos with new styles
  • Editing or extending images by filling masked areas
  • Creating visual content for presentations and projects
  • Building custom image generation pipelines

非目标

  • Real-time video generation
  • Direct image editing without AI assistance
  • Serving models in a production API without additional setup

安装

请先添加 Marketplace

/plugin marketplace add Orchestra-Research/AI-Research-SKILLs
/plugin install AI-Research-SKILLs@ai-research-skills

质量评分

已验证
98 /100
1 day ago 分析

信任信号

最近提交17 days ago
星标8.3k
许可证MIT
状态
查看源代码

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