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Stable Diffusion Image Generation

技能 已验证 活跃

State-of-the-art text-to-image generation with Stable Diffusion models via HuggingFace Diffusers. Use when generating images from text prompts, performing image-to-image translation, inpainting, or building custom diffusion pipelines.

目的

Generate high-quality images from text prompts, perform image transformations, and build custom diffusion workflows using state-of-the-art models.

功能

  • Text-to-Image generation
  • Image-to-Image translation
  • Inpainting and Outpainting
  • ControlNet and LoRA integration
  • Support for various Stable Diffusion models (SDXL, SD 3.0, Flux)

使用场景

  • Generating images from detailed text descriptions
  • Transforming existing images with style or content guidance
  • Filling in masked regions of images contextually
  • Building custom image generation pipelines with advanced controls

非目标

  • API-based generation without local GPU
  • Artistic or highly stylized outputs (Midjourney)
  • Cloud-specific integrations (Imagen, Leonardo.ai)

Trust

  • info:Issues AttentionWith 17 open and 4 closed issues in the last 90 days, the maintainer engagement appears moderate but could be improved for faster issue resolution.

安装

npx skills add davila7/claude-code-templates

通过 npx 运行 Vercel skills CLI(skills.sh)— 需要本地安装 Node.js,以及至少一个兼容 skills 的智能体(Claude Code、Cursor、Codex 等)。前提是仓库遵循 agentskills.io 格式。

质量评分

已验证
95 /100
1 day ago 分析

信任信号

最近提交1 day ago
星标27.2k
许可证MIT
状态
查看源代码

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