Stable Diffusion Image Generation
技能 已验证 活跃State-of-the-art text-to-image generation with Stable Diffusion models via HuggingFace Diffusers. Use when generating images from text prompts, performing image-to-image translation, inpainting, or building custom diffusion pipelines.
Generate high-quality images from text prompts, perform image transformations, and build custom diffusion workflows using state-of-the-art models.
功能
- Text-to-Image generation
- Image-to-Image translation
- Inpainting and Outpainting
- ControlNet and LoRA integration
- Support for various Stable Diffusion models (SDXL, SD 3.0, Flux)
使用场景
- Generating images from detailed text descriptions
- Transforming existing images with style or content guidance
- Filling in masked regions of images contextually
- Building custom image generation pipelines with advanced controls
非目标
- API-based generation without local GPU
- Artistic or highly stylized outputs (Midjourney)
- Cloud-specific integrations (Imagen, Leonardo.ai)
Trust
- info:Issues AttentionWith 17 open and 4 closed issues in the last 90 days, the maintainer engagement appears moderate but could be improved for faster issue resolution.
安装
npx skills add davila7/claude-code-templates通过 npx 运行 Vercel skills CLI(skills.sh)— 需要本地安装 Node.js,以及至少一个兼容 skills 的智能体(Claude Code、Cursor、Codex 等)。前提是仓库遵循 agentskills.io 格式。
质量评分
已验证类似扩展
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